Stable diffusion arts. Jun 5, 2024 · Step 1: Enter a prompt.

Stable Diffusion XL (SDXL) 1. You can modify the prompt below to generate other animals. 4. You will have the opportunity to work on several projects Feb 11, 2024 · To use a VAE in AUTOMATIC1111 GUI, click the Settings tab on the left and click the VAE section. com - Tutorials; Errors. MAT outpainting. Below is an example. Step 3: Remove the triton package in requirements. In this article, you will find a step-by-step guide with detailed Nov 8, 2022 · Software. Luckily, you can use inpainting to fix it. In the inpainting canvas of the img2img tab, draw a mask over the problematic area. Jan 17, 2024 · Step 4: Testing the model (optional) You can also use the second cell of the notebook to test using the model. In general, you should stay away from the two extremes – 1 and 30. Step 6: Convert the output PNG files to video or animated gif. You should now be on the img2img page and Inpaint tab. Jun 5, 2024 · Option 2: Purchase the notebook. Upload an image to the img2img canvas. Living in a vast palace, she lived a privileged life. Step 1. Starting image. Its installation process is no different from any other app. Higher versions have been trained for longer May 25, 2023 · Step-by-step guide. Outpainting complex scenes. See: How Stable Diffusion works. Step 3: Create the animated GIF. Click Install. Dec 21, 2022 · See Software section for set up instructions. Click the ngrok. Gain exposure to media tools like FFmpeg and Davinci Resolve. io in the output under the cell. Using the prompt. Follow these steps to install the Regional Prompter extension in AUTOMATIC1111. Portrait the life of a Roman princess. You will need to sign up to use the model. Method 2: ControlNet img2img. Diffusion is an AI image-generation technique starting with a random image and gradually denoising it to a clear image. It attempts to combine the best of Stable Diffusion and Midjourney: open source, offline, free, and ease-of-use. name is the name of the LoRA model. In the SD VAE dropdown menu, select the VAE file you want to use. You will see the exact keyword applied to two classes of images: (1) a portrait and (2) a scene. . Step 3: Review settings and press generate. Filtering by artists or tags can be done above or by clicking them. We will go with the default setting. Failure example of Stable Diffusion outpainting. Step 1: Generate an initial image. x For the first version 5 model checkpoints are released. All images was generated with the same seed. Switch to img2img tab by clicking img2img. Sep 23, 2023 · tilt-shift photo of {prompt} . Andrew is an experienced engineer with a specialization in Machine Learning and Artificial Intelligence. Some commonly used blocks are Loading a Checkpoint Model, entering a prompt, specifying a sampler, etc. Mar 28, 2023 · The sampler is responsible for carrying out the denoising steps. Step 2: Train a new checkpoint model with Dreambooth. 5 workflow. 30 – Strictly follow the prompt. You don't need to worry. Fooocus has optimized the Stable Diffusion pipeline to deliver excellent images. Mar 19, 2024 · An advantage of using Stable Diffusion is that you have total control of the model. Installing Infinite Zoom on Windows or Mac. 1 means the input image is completely replaced with noise. Nov 26, 2023 · Step 1: Load the text-to-video workflow. Put in a prompt describing your photo. Triggering keyword. It is inspired by the Langevin dynamics formulation of the diffusion process in Physics. Project name. Navigate to the Extension Page. Both modify the U-Net through matrix decomposition, but their approaches differ. 0: (1) Web services, (2) local install and (3) Google Colab. LoRA is the original method. Step 1: Load the FreeU workflow. It is a common setting in image-to-image applications in Stable Diffusion. Upload the photo you want to be cartoonized to the canvas in the img2img sub-tab. All information has been collected with the utmost care, however, mistakes happen. Tweaks for your own artwork. It’s critical to dial up CFG scale to around 11-12. This model allows for image variations and mixing operations as described in Hierarchical Text-Conditional Image Generation with CLIP Latents, and, thanks to its modularity, can be combined with other models such as KARLO. Create stunning AI art for free with SeaArt AI art generator. Next you will need to give a prompt. The model is released as open-source software. Step 1: Select a checkpoint model. Press the big red Apply Settings button on top. 5 sucks at it. Two main ways to train models: (1) Dreambooth and (2) embedding. The steps in this workflow are: Build a base prompt. This approach aims to align with our core values and democratize access, providing users with a variety of options for scalability and quality to best meet their creative needs. This page lists all 1,833 artists that are represented in the Stable Diffusion 1. Keep the denoising strength at 1. Useful Links. Pixel Feb 27, 2023 · Alternative, if you are using the downloaded image, go to img2img tab and select the Inpaint sub-tab. The words it knows are called tokens, which are represented as numbers. IT'S FREE TO USE! · Create artistic masterpieces using text prompts with no need for prior drawing or design skills. 7 – A good balance between following the prompt and freedom. 10 to PATH “) I recommend installing it from the Microsoft store. Number of Epochs. oil painting of zwx in style of van gogh. Step 4: Press Generate. Nov 25, 2023 · The hypernetwork is usually a straightforward neural network: A fully connected linear network with dropout and activation. Apr 27, 2024 · Pixel Art XL is a Stable Diffusion LoRA model available on Civitai that is designed for generating pixel art style images. Option 2: Use the 64-bit Windows installer provided by the Python website. In AUTOMATIC1111 GUI, go to img2img tab and select the img2img sub tab. Fix details with inpainting. Create stunning AI Art in seconds with Stable Diffusion. D. Jun 5, 2024 · Read this post for an overview of the Stable Diffusion 3 model. Use SD3 API with ComfyUI. The suggested animals of this model are pig, bear, chook, monkey, sheep, horse, snake, dragon, bunny, tiger, cow, and rat. Stable Diffusion Version v1. To run Stable Diffusion locally on your PC, download Stable Diffusion from GitHub and the latest checkpoints from HuggingFace. For anime images, it is common to adjust Clip Skip and VAE settings based on the model you use. Stable Diffusion by LMU and stability. Below are a few examples of increasing the CFG scale with the same random seed. They hijack the cross-attention module by inserting two networks to transform the key and query vectors. · Generate AI art in any style (Disney, Pixar, anime, photography, famous artist, anything you can think of) Sep 22, 2022 · You can make NSFW images In Stable Diffusion using Google Colab Pro or Plus. Technically, a positive prompt steers the diffusion toward the images associated with it, while a negative prompt steers the diffusion away from it. Higher versions have been trained for longer Diffusion. It is not a single model but a family of models ranging from 800M to 8B parameters. The prompt should describes both the new style and the content of the original image. 3900+ references. Image Repeats. 5 and 2. The image itself is controlled by the text prompt. We will need the Ultimate SD Upscale and ControlNet extensions for the last method. Jun 5, 2024 · Step 3: Create an API key. Apr 18, 2024 · Fooocus: Stable Diffusion simplified. As each country has their own laws surrounding AI art, it is your responsibility to be compliant. Use the paintbrush tool to create a mask on the face. modyfi : An image editor with AI-powered creative tools for real Sep 27, 2023 · LyCORIS and LoRA models aim to make minor adjustments to a Stable Diffusion model using a small file. You are unauthorized to view this page. Nov 28, 2023 · This is because the face is too small to be generated correctly. Step 2: Enter the text-to-image setting. Dec 14, 2023 · Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. Step 1: Go to DiffusionBee’s download page and download the installer for MacOS – Apple Silicon. Create an image with black text and white background. It saves you time and is great for quickly fixing common issues like garbled faces. Stable Diffusion GUI. Jan 16, 2024 · Option 1: Install from the Microsoft store. This article provides step-by-step guides for creating them in Stable Diffusion. The tags are scraped from Wikidata, a combination of "genres" and "movements". Counterfeit is one of the most popular anime models for Stable Diffusion and has over 200K downloads. It recognizes that many tokens are redundant and can be combined without much consequence. If you put in a word it has not seen before, it will be broken up into 2 or more sub-words until it knows what it is. You need to ask for a specific kind of image. She had everything except freedom and company. Artificial Intelligence (AI) art is currently all the rage, but most AI image generators run in the cloud. 5 base model. Method 3: Dreambooth. Dec 16, 2023 · Windows or Mac. Augmentation level. Step 2. Installing the IP-adapter plus face model. Dec 3, 2023 · When using a negative prompt, a diffusion step is a step towards the positive prompt and away from the negative prompt. Stable Diffusion is a powerful AI image generator. It does not need to be super detailed. Login Username or E-mail Password Remember Me Forgot Password? Click here Apr 13, 2023 · To fix it, first click on Send to inpaint to send the image and the parameters to the inpainting section of the img2img tab. If you're interested in learning more about how to do this using the latter AI-powered tool, I recommend that you take a few minutes to read my article on Midjourney prompts for NFT art. Remove background in Stable Diffusion. New stable diffusion finetune ( Stable unCLIP 2. In my example, I will ask for “photorealistic close-up illustration”. You can create your own model with a unique style if you want. Mar 2, 2023 · In this post, you will see images with diverse styles generated with Stable Diffusion 1. Set Seed to -1 (random), denoising strength to 1, and batch size to 8. Step 4: Choose a seed. Automatic1111 for the Stable Diffusion Web UI. ) Set the Mask Blur to 40. to is an online stable diffusion AI art generator with 8 custom models to choose from. Oct 9, 2023 · Step 1: Install the QR Code Control Model. Step 3: Generate images. On the Settings page, click User Interface on the left panel. Jun 5, 2024 · Stop the instance. Step 3: Enter ControlNet Setting. ai Feb 13, 2024 · SD Upscale is a script that comes with AUTOMATIC1111 that performs upscaling with an upscaler followed by an image-to-image to enhance details. Because of the changes in the language model, prompts that work for SDXL can be a bit different from the v1 models. It's a versatile model that can generate diverse Stable Diffusion Art. Dec 9, 2023 · Welcome to the world of art magic with Stable Diffusion!If you love art and want to make cool pictures using AI, you’re in the perfect spot. They are all generated from simple prompts designed to show the effect of certain keywords. Step 2: Generate SVD video. (If you don’t see this option, you need to update your A1111. weight is the emphasis applied to the LoRA model. Write the prompt and negative prompt in the corresponding input boxes. Soft inpainting seamlessly adds new content that blends with the original image. Sep 8, 2023 · The Stable Diffusion XL (SDXL) model is the latest innovation in the Stable Diffusion technology. Jun 5, 2024 · Soft Inpainting. 98 billion for the v1. Upscale the image. Upload the image by dragging and dropping to the image canvas. ai. Running Stable Diffusion in the cloud (AWS) has many advantages. Stable Diffusion. Mar 21, 2024 · Click the play button on the left to start running. Nov 12, 2023 · 14 Stable Diffusion Prompt Examples for NFT Art. ComfyUI breaks down a workflow into rearrangeable elements so you can easily make your own. Stable Diffusion is a deep learning, text-to-image model released in 2022 based on diffusion techniques. Links that helped me understand the technical side of Stable Diffusion (no affiliation): Automatic111 Wiki; Stable-Diffusion-Art. 6,498 views. Jul 14, 2023 · The Stable Diffusion XL (SDXL) model is the official upgrade to the v1. Stable Diffusion v1. Then run Stable Diffusion in a special python environment using Miniconda. Stable. It also features some great Stable-Diffusion-Art. The courses are truly amazing and easy to follow. The noise predictor then estimates the noise of the image. Nov 7, 2022 · See courses. 5 model. They are LoCon, LoHa, LoKR, and DyLoRA. The first link in the example output below is the ngrok. Sample images. Jul 4, 2023 · Token merging. Both Stable Diffusion XL and Midjourney are excellent AI tools for making NFT art. Jul 6, 2024 · ComfyUI is a node-based GUI for Stable Diffusion. Settings: sd_vae applied. Step 1: Generate an AI image. Step 1: Download an inpainting model. Install Stable Video Diffusion on Windows. Method 4: LoRA. 2. When you visit the ngrok link, it should show a message like below. 3. Nov 29, 2022 · A simple prompt for generating decorative gemstones on fabric. To use the SDXL model, select SDXL Beta in the model menu. The amount of token merging is controlled by the percentage of token merged. 0 means no noise is added to the input image. Extended Artist-style comparison available here. He is passionate about programming, art, photography, and education. You rent the hardware on-demand and only pay for the time you use. First, remove all Python versions you have previously installed. A dmg file should be downloaded. Installing Infinite Zoom on Google Colab. in engineering. Select your Stable Diffusion instance > Instance state > Stop instance. The generative artificial intelligence technology is the premier product of Stability AI and is considered to be a part of the ongoing artificial intelligence boom . mp4 Stable UnCLIP 2. Step 2: Double-click to run the downloaded dmg file in Finder. Choose a model. This page can act as an art reference. Step 1: Clone the repository. art is an open-source plugin for Photoshop (v23. In the AI world, we can expect it to be better. Learn how to create the Women of the World AI art project from start to finish. Step 1: Load the workflow. Jul 22, 2023 · After Detailer (adetailer) is a Stable Diffusion Automatic11111 web-UI extension that automates inpainting and more. It delivers remarkable advancement in image quality. Stable Diffusion XL. Last updated 7/2024. Final adjustment with photo-editing software. Upscale your images, create variations, fix faces, share your art, and more. This beginner’s course introduces you to image-to-image, using checkpoint models correctly, prompt building, A1111 extensions, LoRA, textual inversion, hypernetwork, upscalers, basic ControlNet, and an intermediate end-to-end workflow. Step 2: Load a SDXL model. We will also delve into the exciting process of building your own model using Stable Diffusion. Step 2: Update ComfyUI. Turn on Soft Inpainting by checking the check box next to it. (Alternatively, use Send to Img2img button to send the image to the img2img canvas) Step 3. Generating images with a consistent style is a valuable technique in Stable Diffusion for creative works like logos or book illustrations. 1-768. 15 – Adhere more to prompt. Step 3: Enter ControlNet settings. “close up” and “angled view” did the job. Dec 26, 2023 · Step 2: Select an inpainting model. Use the LoRA with the sunshinemix_sunlightmixPruned model. 4. 3. In the Quicksetting List, add the following. Our design tool is specially created for modern marketers to grow their brands on social media platforms. You will first need to get a fashion model. Jul 5, 2023 · The original image to be stylized. blurry, noisy, deformed, flat, low contrast, unrealistic, oversaturated, underexposed. Jun 12, 2024 · Using LCM-LoRA in AUTOMATIC1111. This model is perfect for generating anime-style images of characters, objects, animals, landscapes, and more. Stable Diffusion XL artists list. See my quick start guide for setting up in Google’s cloud server. 6 (232 ratings) 2,021 students. Say goodbye to tedious art processes and hello to seamless creativity with Stable. Jan 22, 2023 · Close-up illustration. . Step 1: Generate training images with ReActor. Just like the ones you would learn in the introductory course on neural networks. io link to start AUTOMATIC1111. THIS COURSE WILL SHOW YOU HOW TO: · Use the most popular AI Art creating tool: Stable Diffusion. Step 5: Batch img2img with ControlNet. You will get some free credits after signing up. Comparison of all Artists in Stable Diffusion. He has a Ph. Whether you want to create a model that mimics your unique art style or generates images based on your personal photos, we will guide you through the process of bringing your vision to life. Tap into 300K+ models & styles, boost creativity with swift AI tools and engage with the community. Learning rate. By default, Colab notebooks rely on the original Stable Diffusion which comes with NSFW filters. Regardless of the method you use, now you should have setup the GUI like below. With Simplified, you can harness the power of Stable Diffusion as well as DALL-E to create mesmerizing and highly detailed AI images. Nov 16, 2023 · 1 – Mostly ignore your prompt. In other words, the smallest model is a bit smaller than Stable Diffusion 1. Sep 27, 2023 · The workflow is a multiple-step process. Stable Diffusion text-to-image AI art generator. See the following examples of consistent logos created using the technique described in this article. Stable Diffusion 3 combines a diffusion transformer architecture and flow matching. Sep 22, 2023 · The SDXL model is currently available at DreamStudio, the official image generator of Stability AI. Learn how Stable Diffusion predicts noise and how the CFG scale guides the model's prediction. ComfyUI LCM-LoRA SDXL text-to-image workflow. Generating legible text on images has long been challenging for AI image generators. This guide is like your friendly map to creating Jul 31, 2023 · Check out the Quick Start Guide if you are new to Stable Diffusion. Jun 9, 2024 · In text-to-image, you give Stable Diffusion a text prompt, and it returns an image. art. Automatic1111 for the Stable Diffusion web UI. Token merging (ToMe) is a new technique to speed up Stable Diffusion by reducing the number of tokens (in the prompt and negative prompt) that need to be processed. This process is repeated a dozen times. Jan 23, 2024 · Denoising strength determines how much noise is added to an image before the sampling steps. Mentioning an artist in your prompt has a strong influence on generated images. When it is done loading, you will see a link to ngrok. com - Tutorials; Credits. You can construct an image generation workflow by chaining different blocks (called nodes) together. The predicted noise is subtracted from the image. Explore thousands of high-quality Stable Diffusion models, share your AI-generated art, and engage with a vibrant community of creators Dec 19, 2023 · Installing the background removal extension. 5 models. By simply replacing all instances linking to the original script with the script that has no safety filters, you can easily achieve generate NSFW images. You are required to follow the laws of the jurisdication where you live. Jun 5, 2024 · What is the Stable Diffusion 3 model? Stable Diffusion 3 is the latest generation of text-to-image AI models to be released by Stability AI. Step 3: Download models. Patterns are usually generated with a top-down view. It's commonly used for generating artistic images, but can also generate images that look more like photos or sketches. Step 3: Set outpainting parameters. Code from Himuro-Majika's Stable Diffusion image metadata viewer Browser Extension. Table of Contents. Stable Cascade is a quantum leap. Dec 24, 2023 · MP4 video. Jul 4, 2023 · Lonely Palace - Stable Diffusion Art. Gain an understanding of sampling methods and why they are included in the image generation process. Jun 5, 2024 · Step 1: Enter a prompt. If you set the seed to a certain value, you will always get the same random tensor. Additional resources. Step 2: Select the inpainting model. Center an image. Method 5: ControlNet IP-adapter face. Step 2: Enter Img2img settings. (If you use this option, make sure to select “ Add Python to 3. Dreambooth is considered more powerful because it fine-tunes the weight of the whole model. Method 2: Generate a QR code with the tile resample model in image-to-image. Step 3: Download and load the LoRA. Step 2: Remove background. You control this tensor by setting the seed of the random number generator. A BIG THANK YOU TO. Jan 4, 2024 · The CLIP model Stable Diffusion automatically converts the prompt into tokens, a numerical representation of words it knows. Fix the subject. Navigate to Img2img page. Nov 5, 2023 · Stable Diffusion Software. Step 2: Create a virtual environment. Fooocus is a free and open-source AI image generator based on Stable Diffusion. Aug 16, 2023 · Tips for using ReActor. Click the Available tab. 0. 0+) that allows you to use Stable Diffusion (with Automatic1111 as a backend) to accelerate your art workflow. with my newly trained model, I am happy with what I got: Images from dreambooth model. Text rendering. Stable Diffusion XL is an improvement. You can use this GUI on Windows, Mac, or Google Colab. It is infinitely customizable and always unique. Use FreeU in ComfyUI. Adding the LCM sampler with AnimateDiff extension. Extensions shape our workflow and make Stable Diffusion even more powerful. Step 3: Define system-wide API key (optional but recommended) Step 4: Load and run the workflow. Click Load from: button. “bokeh” add to close up view nicely. You can use this website or any photo editing tool. 6 billion, compared with 0. It's created by researchers and engineers from CompVis, Laion, and Stability AI. In AUTOMATIC1111 GUI, select the Inpunk Diffusion model in the Stable Diffusion checkpoint dropdown menu. Summary. Step 2: Come up with a good triggering keyword. Dec 21, 2022 · Images from celebrity and commercial artists were also suppressed. Upload the image to the img2img canvas. Motion bucket id. Step 1: Convert the mp4 video to png files. Using the IP-adapter plus face model. Its community-developed extensions make it stand out, enhancing its functionality and ease of use. I wanted to create a oblique view to make it more interesting. The best tools to generate this kind of image are ControlNet and ADetailer. Step 4: Generate images. When you are done, stop the instance to avoid extra charges. Start AUTOMATIC1111 Web-UI normally. You can click on an image to enlarge it. Step 3: Generate image. Prompt: oil painting of zwx in style of van gogh. Nov 22, 2023 · To add a LoRA with weight in AUTOMATIC1111 Stable Diffusion WebUI, use the following syntax in the prompt or the negative prompt: <lora: name: weight>. Step 3: Review the training settings. Selecting the SDXL Beta model in DreamStudio. Step 2: Select a checkpoint model. List of artists supported by Stable May 12, 2023 · Software. In this post, I will go through the workflow step-by-step. We will use AUTOMATIC1111 Stable Diffusion GUI to perform upscaling. In the second part, I will compare images generated with Stable Diffusion 1. It is a much larger model. Apr 29, 2023 · The Chinese Zodiac LoRA generates cute animals in a cartoon style. We will use Stable Diffusion AI and AUTOMATIC1111 GUI. SVD settings. Stable Diffusion is an advanced AI text-to-image synthesis algorithm that can generate very coherent images based on a text prompt. Feb 16, 2023 · Key Takeaways. 1. Updated July 4, 2023By Andrew Workflow Tagged Historical 2 Comments on Lonely Palace. 0 is Stable Diffusion's next-generation model. Find the extension “Regional Prompter”. Return mask. Refinement prompt and generate image with good composition. Convert to landscape size. Alpha matting. co, and install them. Stable Diffusion 1. May 15, 2024 · DiffusionBee is one of the easiest ways to run Stable Diffusion on Mac. In this post, you will learn how it works, how to use it, and some common use cases. Set the image size to 768 x 512 pixels. Can’t ask for anything better! Oct 5, 2023 · Simplified is one of the best AI art and image generator in 2024. FPS. English. Open the Amazon EC2 console. Dec 10, 2023 · Pixelz AI Art Generator: Enables the creation of art from text using algorithms such as Stable Diffusion and CLIP Guided Diffusion. Created by Chris White. Only artist name was changed in prompts. demo. Fix defects with inpainting. There are three important techniques to tease out high-quality prompts for Stable Diffusion from ChatGPT: Specify image style. 1, Hugging Face) at 768x768 resolution, based on SD2. Background threshold. Step 1: Prepare training images. Open AUTOMATIC1111 WebUI and navigate to the txt2img page. 5 (1B), and the largest model is a bit bigger Feb 22, 2024 · The Stable Diffusion 3 suite of models currently ranges from 800M to 8B parameters. The good news is Stable Diffusion 3’s text generation is at the next level. The total number of parameters of the SDXL model is 6. Stable Diffusion Level 2. 4 Model, ordered by the frequency of their representation. Step 2: Enter a prompt and a negative prompt. io link. Notes for ControlNet m2m script. Apr 13, 2024 · Install SVD model. Foreground threshold. It aims to produce consistent pixel sizes and more “pixel perfect” outputs compared to standard Stable Diffusion models. Step 4: Run the workflow. Step 1: Install ComfyUI Manager. Note that the diffusion in Stable Diffusion happens in latent space, not images. Advanced options. Use the SVD model to generate a video. You can also combine it with LORA models to be more versatile and generate unique artwork. Become a Stable Diffusion Pro step-by-step. It can be used with Stable Diffusion XL (SDXL) models to generate pixel art style images. LyCORIS is a collection of LoRA-like methods. It can be different from the filename. Step 1: Add a waterfall. In this article, I will cover 3 ways to run Stable diffusion 2. Select a starting image. Save as a PNG file. Mar 10, 2024 · Apr 29, 2023. Feb 18, 2024 · Must-have AUTOMATIC1111 extensions. Stable Diffusion generates a random tensor in the latent space. Restart the web-ui. The value of denoising strength ranges from 0 to 1. It is similar to a keyword weight. 3 – Be more creative. A simple first example. To produce an image, Stable Diffusion first generates a completely random image in the latent space. Prompt: Ultimate Stable Diffusion AI art Course (beginner to pro) Learn prompt, control character pose and lighting, train your own model, ChatGPT and more with Stable Diffusion! Bestseller. Step 3: Using the model. It is convenient to enable them in Quick Settings. Step 2: Enable FreeU. Jun 5, 2024 · Stable Diffusion AI is particularly suited for generating seemingly innocent images with hidden words. Lazyload Script from Verlok, Webfont is Google's Roboto, SVG Icons from Ionicons, MattiasW's ExifReader. Software. Step 4: Enable the outpainting script. Step 2: Animate the water fall and the clouds. Diffusion. Step 2: Install the SAI API node. selective focus, miniature effect, blurred background, highly detailed, vibrant, perspective control. You should see the message. Lonely Palace. Click the Send to Inpaint icon below the image to send the image to img2img > inpainting. Parameters. fb ri zb hr qz uu aq ho gj dl